Latest Models

Effortlessly merge audio and video with our intuitive Video Audio Merge model. Create stunning multimedia content with precise timing, fade effects, and customizable audio options. Perfect for content creators, filmmakers, and marketers.

Runway Gen-3 AlphaTurbo is a cutting-edge AI tool that transforms static images into dynamic videos with exceptional fidelity and motion

MusicGen: Transform text into music with AI. Create unique, high-quality audio from simple descriptions. Experience the future of music generation with this innovative AI model.

With Video Captioner create accurate, customizable subtitles for your videos effortlessly.

Restore characters' faces to their original glory with Face Detailer. Enhance facial details, eliminate distortion, and upscale images for stunning results.

Flux 1.1 Pro is a cutting-edge image generation tool offering exceptional speed, quality, and customization. Ideal for digital artists, designers, and content creators.

MyShell's Voice Cloning and Text to Speech - Transform your audio content with realistic, personalized voices. Experience high-quality, efficient, and cost-effective audio synthesis.

Revolutionize your video editing with the Video Stitch Model. Seamlessly stitch clips, add captivating audio, and create professional-looking videos in minutes.

Flux is a 12 billion parameter rectified flow transformer capable of generating images from text descriptions

Ideogram Text to Image: Turn your ideas into stunning visuals instantly with this powerful AI tool. Create captivating designs, realistic images, and more. Perfect for artists, designers, and anyone seeking creative inspiration.

OpenVoice is a versatile voice cloning model that supports multiple languages and offers precise tone replication, flexible style control, and zero-shot cross-lingual capabilities

CogVideoX image-to-video is a cutting-edge AI model that converts static images into dynamic, high-quality videos. Perfect for content creation, animation, and education, it offers high-resolution output, efficient inference, and versatile precision. Transform your images into engaging videos with CogVideoX

Expression Editor uses reference images to accurately generate new images with desired expressions. Perfect for digital art, memes, and marketing.

ESRGAN Video Upscaler: Experience sharper, clearer videos with ESRGAN. This AI-powered video upscaler boosts resolution and reduces artifacts, making your content look its best.

Create images of a given character in different poses

Create consistent characters in any pose with AI

Luma Video (Text to Video) is an advanced AI model that turns text prompts into captivating videos. Designed for creators and marketers, it offers high-resolution outputs, rapid processing, and cinematic quality, making video production accessible and efficient.

With Luma's Dream Machine, transform your static images into dynamic videos. It offers high-fidelity video generation, rapid processing, and cinematic quality, enabling users to enhance their content creation process effortlessly.

Flux PuLID: Customize AI-generated images with your unique identity. Seamlessly integrate faces into text-to-image models for realistic and customizable results. High fidelity, tuning-free customization, and versatile editing options.

Flux IP Adapter is a cutting-edge AI model that lets you to create stunning, customized images. With its advanced style adaptation capabilities, Flux IP Adapter lets you seamlessly blend different artistic styles into your creations.

Flux Inpainting is a powerful image editing tool designed to effortlessly edit and enhance your images. It's perfect for tasks like removing unwanted objects, restoring damaged photos, and creating artistic effects.

Flux ControlNets is a collection of models that gives you precise control over image generation. By integrating ControlNet with Flux.1, these models enable you to create highly detailed and customized images with unprecedented accuracy.

o1-mini by OpenAI provides high-performance reasoning and coding capabilities. Ideal for developers and businesses seeking advanced AI without the high costs.

o1-preview by OpenAI, is a powerful AI model that can tackle complex problems with exceptional accuracy and efficiency. Ideal for researchers, developers, and scientists seeking advanced AI capabilities.

Elevate your visuals withText Overlay Model. Easily add customized text to any image, perfect for social media, marketing, and blogs. Enjoy precise positioning, advanced styling, and seamless integration.

CogVideo is a groundbreaking AI model that turns text into high-quality videos. Create realistic scenes, animations, and more with ease. Ideal for content creators, educators, and businesses.

Fast Flux.1 Schnell by Segmind is an optimized text-to-image model designed for developers needing faster image generation. It offers high efficiency without compromising quality. Perfect for startups and engineers seeking quick, resource-efficient AI models.

Create consistent characters in any pose with AI

Flux Realism Lora with upscale, developed by XLabs AI is a cutting-edge model designed to generate realistic images from textual descriptions.

SAM v2, the next-gen segmentation model from Meta AI, revolutionizes computer vision. Building on SAM's success, it excels at accurately segmenting objects in images, offering robust and efficient solutions for various visual contexts.

SAM v2 Video by Meta AI, allows promptable segmentation of objects in videos.

Create images of a given character in different poses

Flux Image-To-Image model by Black Forest Labs is an advanced deep learning tool designed for transforming images based on specific textual prompts.

Easy Animate is a state-of-the-art image to animation model to convert static images into dynamic animations with remarkable accuracy and fluidity.

Flux Dev is a 12 billion parameter rectified flow transformer capable of generating images from text descriptions

Flux Schnell  is a state-of-the-art text-to-image generation model engineered for speed and efficiency.

Flux Pro is a state-of-the-art image generation with top of the line prompt following, visual quality, image detail and output diversity.

Text-embedding-3-small is a compact and efficient model developed for generating high-quality text embeddings. These embeddings are numerical representations of text data, enabling a variety of natural language processing (NLP) tasks such as semantic search, clustering, and text classification

Text-embedding-3-large is a robust language model by OpenAI designed for generating high-dimensional text embeddings for a wide range of natural language processing (NLP) tasks including semantic search, text clustering, and classification.

Real Dream Pony V9 is an advanced image generation model based on the Stable Diffusion XL (SDXL) architecture, excelling in photorealism.

AI Product Photo Editor leverages advanced image-based ML techniques to generate high-quality product visuals using text prompts, product images, and background images.

RealDream is a sophisticated image generation model utilizing SDXL Lightning architecture. It creates incredibly realistic images from textual prompts. With the ability to excellently generate human portraits from the user's descriptive text.

Meta developed and released the Meta Llama 3 family of large language models (LLMs), a collection of pretrained and instruction tuned generative text models in 8 and 70B sizes. The Llama 3 instruction tuned models are optimized for dialogue use cases and outperform many of the available open source chat models on common industry benchmarks.

Meta developed and released the Meta Llama 3 family of large language models (LLMs), a collection of pretrained and instruction tuned generative text models in 8 and 70B sizes. The Llama 3 instruction tuned models are optimized for dialogue use cases and outperform many of the available open source chat models on common industry benchmarks.

Meta developed and released the Meta Llama 3 family of large language models (LLMs), a collection of pretrained and instruction tuned generative text models in 8 and 70B sizes. The Llama 3 instruction tuned models are optimized for dialogue use cases and outperform many of the available open source chat models on common industry benchmarks.

Stable Diffusion 3 Medium image-to-image is a cutting-edge AI tool that uses advanced image-to-image technology to transform one image into another.

SD3 Medium Tile ControlNet is a large generative image model designed for generating detailed images based on textual prompts and tile-based input images.

Stable Diffusion 3 (SD3) Medium Canny ControlNet uses Canny edge detection to provide fine-grained control over the generated outputs.

Stable Diffusion 3 (SD3) Pose ControlNet is a large generative image model tailored for generating images based on text prompts while using pose information as guidance.

Motion Control SVD is an innovative deep learning framework that breathes life into static images. By intelligently managing both camera and object motion, it empowers creators to achieve precise animation effects.

Experience the magic of Live Portrait’s Video-to-Video Model! Transform your static images into dynamic videos seamlessly.

Superimpose V2 elevates image editing! Seamlessly layer images with background removal, precise positioning, and flexible resizing options. Explore 14 blending modes to create stunning effects

Video Faceswap is a powerful tool for creators, filmmakers, and meme enthusiasts. With this innovative technology, you can effortlessly replace faces in videos

Largest completely open sourced flow-based generation model that is capable of text-to-image generation

Live Portrait animates static images using a reference driving video through implicit key point based framework, bringing a portrait to life with realistic expressions and movements. It identifies key points on the face (think eyes, nose, mouth) and manipulates them to create expressions and movements.

ElevenLabs Dubbing uses AI to translate your audio into multiple languages. Easily create multilingual versions of your content without studios or voice actors for each language

Claude 3 Haiku, the fastest and most cost-effective model LLM from Anthropic, delivers instant responses and image analysis. Build interactive AI experiences that mimic human conversation. Perfect for various applications, from research to enterprise

Claude 3 Opus is an LLM pushing the limits of language understanding. It excels at complex tasks, generates human-quality text, and remembers vast amounts of information.

Gemini 1.5 Pro represents a significant leap in large language model technology, offering exceptional understanding and performance across different modalities and contexts.

Gemini 1.5 Flash is a game-changer for developers and enterprises seeking a speedy and cost-effective large language model with exceptional long-context understanding.

Claude 3.5 Sonnet represents a significant advancement in AI language models, combining speed, accuracy, and visual reasoning capabilities. It excels at understanding and completing requests thoughtfully, and does so much faster than previous versions. Additionally, it boasts a stronger vision model, allowing it to analyze visual data like charts and images with exceptional accuracy.

Kolors is a cutting-edge text-to-image model that bridges language and visual art. Transform your textual ideas into photorealistic images with semantic precision.

Playground V2.5 is a diffusion-based text-to-image generative model, designed to create highly aesthetic images based on textual prompts.

Superimpose model lets you to create captivating visuals by seamlessly overlaying one image on top of another. It streamlines your image layering process, allowing you to bring your creative vision to life effortlessly.

SDXL Img2Img is used for text-guided image-to-image translation. This model uses the weights from Stable Diffusion to generate new images from an input image using StableDiffusionImg2ImgPipeline from diffusers

SDXL ControlNet gives unprecedented control over text-to-image generation. SDXL ControlNet models Introduces the concept of conditioning inputs, which provide additional information to guide the image generation process

Story Diffusion turns your written narratives into stunning image sequences.

Eleven Labs' Sound Generation API provides a robust development tool for programmatically generating audio content using artificial intelligence. This API empowers developers and creators to integrate sound generation functionalities into their applications and workflows.

Eleven Labs Speech-to-Speech offers AI-powered voice conversion for content creators, media professionals, and anyone seeking to modify or translate audio speech.

Eleven Labs Text-to-Speech (TTS) harnesses the power of deep learning to create realistic and engaging synthetic speech from written text.

Omni-Zero: A diffusion pipeline for zero-shot stylized portrait creation.

LLaVa translates images into text descriptions & captions.

LLaVA 13B is a Vision-language model which allows both image and text as inputs.

Create videos from illustrated input images

V-Express lets you create portrait videos from single images.

Audio-based Lip Synchronization for Talking Head Video

Hallo lets you create portrait videos from single images.

Prompts to auto-magically relight your images.

Automatic Mask Generator is a powerful tool that automates the creation of precise masks for inpainting

LaMA Object Removal- AI Magic Eraser

Real-Time Open-Vocabulary Object Detection

Background Eraser helps in flawless background removal with exceptional accuracy.

High resolution creative image Upscaler and Enhancer. A free Magnific alternative.

Create images of a given character in different poses

Best-in-class clothing virtual try on in the wild

Stable Diffusion is a type of latent diffusion model that can generate images from text. It was created by a team of researchers and engineers from CompVis, Stability AI, and LAION. Stable Diffusion v2 is a specific version of the model architecture. It utilizes a downsampling-factor 8 autoencoder with an 865M UNet and OpenCLIP ViT-H/14 text encoder for the diffusion model. When using the SD 2-v model, it produces 768x768 px images. It uses the penultimate text embeddings from a CLIP ViT-H/14 text encoder to condition the generation process.

Fooocus enables high-quality image generation effortlessly, combining the best of Stable Diffusion and Midjourney.

Style & Composition Transfer with Stable Diffusion IP Adapter

Turn any image of a face into artwork using Stable Diffusion Controlnet and IPAdapter

Monster Labs QrCode ControlNet on top of SD Realistic Vision v5.1

Novel tuning-free ID customization method for text-to-image generation.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

GPT-4 outperforms both previous large language models and as of 2023, most state-of-the-art systems (which often have benchmark-specific training or hand-engineering). On the MMLU benchmark, an English-language suite of multiple-choice questions covering 57 subjects, GPT-4 not only outperforms existing models by a considerable margin in English, but also demonstrates strong performance in other languages. Currently points to gpt-4-turbo-2024-04-09.

GPT-4o (“o” for “omni”) is our most advanced model. It is multimodal (accepting text or image inputs and outputting text), and it has the same high intelligence as GPT-4 Turbo but is much more efficient—it generates text 2x faster and is 50% cheaper. Additionally, GPT-4o has the best vision and performance across non-English languages of any of our models. GPT-4o is available in the OpenAI API to paying customers.

GPT-4 outperforms both previous large language models and as of 2023, most state-of-the-art systems (which often have benchmark-specific training or hand-engineering). On the MMLU benchmark, an English-language suite of multiple-choice questions covering 57 subjects, GPT-4 not only outperforms existing models by a considerable margin in English, but also demonstrates strong performance in other languages.

Mistral MoE 8x7B Instruct v0.1 model with Sparse Mixture of Experts. Fine tuned for instruction following.

Mistral MoE 8x22B Instruct v0.1 model with Sparse Mixture of Experts. Fine tuned for instruction following.

Faster version of PuLID, a novel tuning-free face customization method for text-to-image generation

This model is capable of editing clothing in an image using a premier clothing segmentation algorithm.

Turn a face into 3D, emoji, pixel art, video game, claymation or toy

Turn a face into a sticker

Meta developed and released the Meta Llama 3 family of large language models (LLMs), a collection of pretrained and instruction tuned generative text models in 8 and 70B sizes. The Llama 3 instruction tuned models are optimized for dialogue use cases and outperform many of the available open source chat models on common industry benchmarks.

Transfer a material from an image to a subject

Meta developed and released the Meta Llama 3 family of large language models (LLMs), a collection of pretrained and instruction tuned generative text models in 8 and 70B sizes. The Llama 3 instruction tuned models are optimized for dialogue use cases and outperform many of the available open source chat models on common industry benchmarks.

Take a picture/gif and replace the face in it with a face of your choice. You only need one image of the desired face. No dataset, no training

InstantID aims to generate customized images with various poses or styles from only a single reference ID image while ensuring high fidelity

This model removes the background image from any image

NewReality Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

DreamShaper Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Colossus Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Samaritan Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Realism Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

ProtoVision Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

NightVis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

WildCard Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Dynavis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Juggernaut Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Realvis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Outfitting Fusion based Latent Diffusion for Controllable Virtual Try-on

This model is capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask

Fooocus Inpainting is a powerful image generation model that allows you to selectively edit and enhance images.

Fooocus Outpainting transforms ordinary images into extraordinary works of art by seamlessly expanding their boundaries.

InstantID aims to generate customized images with various poses or styles from only a single reference ID image while ensuring high fidelity

Samaritan 3D XL leverages the robust capabilities of the SDXL framework, ensuring high-quality, detailed 3D character renderings.

Takes image as input and returns a video.

The Segmind-Vega Model is a distilled version of the Stable Diffusion XL (SDXL), offering a remarkable 70% reduction in size and an impressive 100% speedup while retaining high-quality text-to-image generation capabilities.

Segmind-VegaRT a distilled consistency adapter for Segmind-Vega that allows to reduce the number of inference steps to only between 2 - 8 steps.

IP Adapter XL Openpose is built on the SDXL framework. This model integrates the IP Adapter and Openpose preprocessor to offer unparalleled control and guidance in creating context-rich images.

IP Adpater XL Canny is built on the SDXL framework. This model integrates the IP Adapter and Canny edge preprocessor to offer unparalleled control and guidance in creating context-rich images.

IP Adapter Depth XL is built on the SDXL framework. This model integrates the IP Adapter and Depth preprocessor to offer unparalleled control and guidance in creating context-rich images.

This model is capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask

This model uses SSD-1B to generate images by passing a text prompt and an initial image to condition the generation

This model leverages SDXL to generate the images with ControlNet conditioned on Human Pose Estimation.

This model leverages SSD-1B to generate the images with ControlNet conditioned on Depth Estimation

This model leverages SSD-1B to generate the images with ControlNet conditioned on Canny Images

The Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of the Stable Diffusion XL (SDXL), offering a 60% speedup while maintaining high-quality text-to-image generation capabilities. It has been trained on diverse datasets, including Grit and Midjourney scrape data, to enhance its ability to create a wide range of visual content based on textual prompts.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

Stable Diffusion is a type of latent diffusion model that can generate images from text. It was created by a team of researchers and engineers from CompVis, Stability AI, and LAION. Stable Diffusion v2 is a specific version of the model architecture. It utilizes a downsampling-factor 8 autoencoder with an 865M UNet and OpenCLIP ViT-H/14 text encoder for the diffusion model. When using the SD 2-v model, it produces 768x768 px images. It uses the penultimate text embeddings from a CLIP ViT-H/14 text encoder to condition the generation process.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

Create beautifully designed words using Segmind’s word to image for your marketing purposes

Convert Text into Images with the latest distilled stable diffusion model

Stable Diffusion Inpainting is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask

This model uses diffusion-denoising mechanism as first proposed by SDEdit, Stable Diffusion is used for text-guided image-to-image translation. This model uses the weights from Stable Diffusion to generate new images from an input image using StableDiffusionImg2ImgPipeline from diffusers

Create realistic portrait images using the finetined Segmind Tiny SD model. Segmind Tiny SD (Portrait) Serverless APIs, Segmind offers fastest deployment for Tiny-Stable-Diffusion inferences

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images

This model corresponds to the Stable Diffusion RPG checkpoint for detailed images at the cost of a super detailed prompt

This model corresponds to the Stable Diffusion Reliberate checkpoint for detailed images at the cost of a super detailed prompt

This model corresponds to the Stable Diffusion Realistic Vision checkpoint for detailed images at the cost of a super detailed prompt

The most versatile photorealistic model that blends various models to achieve the amazing realistic images

This model corresponds to the Stable Diffusion Paragon checkpoint for detailed images at the cost of a super detailed prompt

Stable Diffusion Outpainting can extend any image in any direction

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

This model corresponds to the Stable Diffusion Epic Realism checkpoint for detailed images at the cost of a super detailed prompt

This model corresponds to the Stable Diffusion Edge of Realism checkpoint for detailed images at the cost of a super detailed prompt

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

Dreamshaper excels in delivering high-quality, detailed images. It is fine-tuned to understand and interpret a diverse range of artistic styles and subjects.

The most versatile photorealistic model that blends various models to achieve the amazing realistic space themed images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

This model corresponds to the Stable Diffusion Colorful checkpoint for detailed images at the cost of a super detailed prompt

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

Create beautiful and creative QR codes for your marketing campaigns.

Convert text to images with the distilled stable diffusion model by Segmind, Small-SD. Segmind Small SD Serverless APIs, Segmind offers fastest deployment for Small-Stable-Diffusion inferences.

Kandinsky inherits best practices from Dall-E 2 and Latent diffusion, while introducing some new ideas.

This model corresponds to the ControlNet conditioned on Soft Edge.

This model corresponds to the ControlNet conditioned on Scribble images.

This model corresponds to the ControlNet conditioned on Depth estimation.

This model corresponds to the ControlNet conditioned on Canny edges.

CodeFormer is a robust face restoration algorithm for old photos or AI-generated faces.

The Segment Anything Model (SAM) produces high quality object masks from input prompts such as points or boxes, and it can be used to generate masks for all objects in an image.

Take a picture/gif and replace the face in it with a face of your choice. You only need one image of the desired face. No dataset, no training

This model corresponds to the Stable Diffusion Revanimated checkpoint for detailed images at the cost of a super detailed prompt

This model removes the background image from any image

AI-Powered Image Super-Resolution, upscaling and Image enhancement producing stunning, high-quality results using artificial intelligence

This model corresponds to the ControlNet conditioned on Human Pose Estimation.