Latest Models

Superimpose V2 elevates image editing! Seamlessly layer images with background removal, precise positioning, and flexible resizing options. Explore 14 blending modes to create stunning effects

Video Faceswap is a powerful tool for creators, filmmakers, and meme enthusiasts. With this innovative technology, you can effortlessly replace faces in videos

Largest completely open sourced flow-based generation model that is capable of text-to-image generation

Live Portrait animates static images using a reference driving video through implicit key point based framework, bringing a portrait to life with realistic expressions and movements. It identifies key points on the face (think eyes, nose, mouth) and manipulates them to create expressions and movements.

ElevenLabs Dubbing uses AI to translate your audio into multiple languages. Easily create multilingual versions of your content without studios or voice actors for each language

Claude 3 Haiku, the fastest and most cost-effective model LLM from Anthropic, delivers instant responses and image analysis. Build interactive AI experiences that mimic human conversation. Perfect for various applications, from research to enterprise

Claude 3 Opus is an LLM pushing the limits of language understanding. It excels at complex tasks, generates human-quality text, and remembers vast amounts of information.

Gemini 1.5 Pro represents a significant leap in large language model technology, offering exceptional understanding and performance across different modalities and contexts.

Gemini 1.5 Flash is a game-changer for developers and enterprises seeking a speedy and cost-effective large language model with exceptional long-context understanding.

Claude 3.5 Sonnet represents a significant advancement in AI language models, combining speed, accuracy, and visual reasoning capabilities. It excels at understanding and completing requests thoughtfully, and does so much faster than previous versions. Additionally, it boasts a stronger vision model, allowing it to analyze visual data like charts and images with exceptional accuracy.

Kolors is a cutting-edge text-to-image model that bridges language and visual art. Transform your textual ideas into photorealistic images with semantic precision.

Playground V2.5 is a diffusion-based text-to-image generative model, designed to create highly aesthetic images based on textual prompts.

Superimpose model lets you to create captivating visuals by seamlessly overlaying one image on top of another. It streamlines your image layering process, allowing you to bring your creative vision to life effortlessly.

SDXL Img2Img is used for text-guided image-to-image translation. This model uses the weights from Stable Diffusion to generate new images from an input image using StableDiffusionImg2ImgPipeline from diffusers

SDXL ControlNet gives unprecedented control over text-to-image generation. SDXL ControlNet models Introduces the concept of conditioning inputs, which provide additional information to guide the image generation process

Story Diffusion turns your written narratives into stunning image sequences.

Eleven Labs' Sound Generation API provides a robust development tool for programmatically generating audio content using artificial intelligence. This API empowers developers and creators to integrate sound generation functionalities into their applications and workflows.

Eleven Labs Speech-to-Speech offers AI-powered voice conversion for content creators, media professionals, and anyone seeking to modify or translate audio speech.

Eleven Labs Text-to-Speech (TTS) harnesses the power of deep learning to create realistic and engaging synthetic speech from written text.

Omni-Zero: A diffusion pipeline for zero-shot stylized portrait creation.

LLaVa translates images into text descriptions & captions.

LLaVA 13B is a Vision-language model which allows both image and text as inputs.

Create videos from illustrated input images

V-Express lets you create portrait videos from single images.

Audio-based Lip Synchronization for Talking Head Video

Hallo lets you create portrait videos from single images.

Prompts to auto-magically relight your images.

Automatic Mask Generator is a powerful tool that automates the creation of precise masks for inpainting

LaMA Object Removal- AI Magic Eraser

Real-Time Open-Vocabulary Object Detection

Background Eraser helps in flawless background removal with exceptional accuracy.

High resolution creative image Upscaler and Enhancer. A free Magnific alternative.

Create images of a given character in different poses

Best-in-class clothing virtual try on in the wild

Stable Diffusion is a type of latent diffusion model that can generate images from text. It was created by a team of researchers and engineers from CompVis, Stability AI, and LAION. Stable Diffusion v2 is a specific version of the model architecture. It utilizes a downsampling-factor 8 autoencoder with an 865M UNet and OpenCLIP ViT-H/14 text encoder for the diffusion model. When using the SD 2-v model, it produces 768x768 px images. It uses the penultimate text embeddings from a CLIP ViT-H/14 text encoder to condition the generation process.

Fooocus enables high-quality image generation effortlessly, combining the best of Stable Diffusion and Midjourney.

Style & Composition Transfer with Stable Diffusion IP Adapter

Turn any image of a face into artwork using Stable Diffusion Controlnet and IPAdapter

Monster Labs QrCode ControlNet on top of SD Realistic Vision v5.1

Novel tuning-free ID customization method for text-to-image generation.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

GPT-4 outperforms both previous large language models and as of 2023, most state-of-the-art systems (which often have benchmark-specific training or hand-engineering). On the MMLU benchmark, an English-language suite of multiple-choice questions covering 57 subjects, GPT-4 not only outperforms existing models by a considerable margin in English, but also demonstrates strong performance in other languages. Currently points to gpt-4-turbo-2024-04-09.

GPT-4o (“o” for “omni”) is our most advanced model. It is multimodal (accepting text or image inputs and outputting text), and it has the same high intelligence as GPT-4 Turbo but is much more efficient—it generates text 2x faster and is 50% cheaper. Additionally, GPT-4o has the best vision and performance across non-English languages of any of our models. GPT-4o is available in the OpenAI API to paying customers.

GPT-4 outperforms both previous large language models and as of 2023, most state-of-the-art systems (which often have benchmark-specific training or hand-engineering). On the MMLU benchmark, an English-language suite of multiple-choice questions covering 57 subjects, GPT-4 not only outperforms existing models by a considerable margin in English, but also demonstrates strong performance in other languages.

Mistral MoE 8x7B Instruct v0.1 model with Sparse Mixture of Experts. Fine tuned for instruction following.

Mistral MoE 8x22B Instruct v0.1 model with Sparse Mixture of Experts. Fine tuned for instruction following.

Faster version of PuLID, a novel tuning-free face customization method for text-to-image generation

This model is capable of editing clothing in an image using a premier clothing segmentation algorithm.

Turn a face into 3D, emoji, pixel art, video game, claymation or toy

Turn a face into a sticker

Meta developed and released the Meta Llama 3 family of large language models (LLMs), a collection of pretrained and instruction tuned generative text models in 8 and 70B sizes. The Llama 3 instruction tuned models are optimized for dialogue use cases and outperform many of the available open source chat models on common industry benchmarks.

Transfer a material from an image to a subject

Meta developed and released the Meta Llama 3 family of large language models (LLMs), a collection of pretrained and instruction tuned generative text models in 8 and 70B sizes. The Llama 3 instruction tuned models are optimized for dialogue use cases and outperform many of the available open source chat models on common industry benchmarks.

Take a picture/gif and replace the face in it with a face of your choice. You only need one image of the desired face. No dataset, no training

InstantID aims to generate customized images with various poses or styles from only a single reference ID image while ensuring high fidelity

This model removes the background image from any image

NewReality Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

DreamShaper Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Colossus Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Samaritan Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Realism Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

ProtoVision Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

NightVis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

WildCard Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Dynavis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Juggernaut Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Realvis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.

Outfitting Fusion based Latent Diffusion for Controllable Virtual Try-on

This model is capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask

InstantID aims to generate customized images with various poses or styles from only a single reference ID image while ensuring high fidelity

Samaritan 3D XL leverages the robust capabilities of the SDXL framework, ensuring high-quality, detailed 3D character renderings.

Takes image as input and returns a video.

The Segmind-Vega Model is a distilled version of the Stable Diffusion XL (SDXL), offering a remarkable 70% reduction in size and an impressive 100% speedup while retaining high-quality text-to-image generation capabilities.

Segmind-VegaRT a distilled consistency adapter for Segmind-Vega that allows to reduce the number of inference steps to only between 2 - 8 steps.

IP Adapter XL Openpose is built on the SDXL framework. This model integrates the IP Adapter and Openpose preprocessor to offer unparalleled control and guidance in creating context-rich images.

IP Adpater XL Canny is built on the SDXL framework. This model integrates the IP Adapter and Canny edge preprocessor to offer unparalleled control and guidance in creating context-rich images.

IP Adapter Depth XL is built on the SDXL framework. This model integrates the IP Adapter and Depth preprocessor to offer unparalleled control and guidance in creating context-rich images.

This model is capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask

This model uses SSD-1B to generate images by passing a text prompt and an initial image to condition the generation

This model leverages SDXL to generate the images with ControlNet conditioned on Human Pose Estimation.

This model leverages SSD-1B to generate the images with ControlNet conditioned on Depth Estimation

This model leverages SSD-1B to generate the images with ControlNet conditioned on Canny Images

The Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of the Stable Diffusion XL (SDXL), offering a 60% speedup while maintaining high-quality text-to-image generation capabilities. It has been trained on diverse datasets, including Grit and Midjourney scrape data, to enhance its ability to create a wide range of visual content based on textual prompts.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

Stable Diffusion is a type of latent diffusion model that can generate images from text. It was created by a team of researchers and engineers from CompVis, Stability AI, and LAION. Stable Diffusion v2 is a specific version of the model architecture. It utilizes a downsampling-factor 8 autoencoder with an 865M UNet and OpenCLIP ViT-H/14 text encoder for the diffusion model. When using the SD 2-v model, it produces 768x768 px images. It uses the penultimate text embeddings from a CLIP ViT-H/14 text encoder to condition the generation process.

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.

Create beautifully designed words using Segmind’s word to image for your marketing purposes

Convert Text into Images with the latest distilled stable diffusion model

Stable Diffusion Inpainting is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask

This model uses diffusion-denoising mechanism as first proposed by SDEdit, Stable Diffusion is used for text-guided image-to-image translation. This model uses the weights from Stable Diffusion to generate new images from an input image using StableDiffusionImg2ImgPipeline from diffusers

Create realistic portrait images using the finetined Segmind Tiny SD model. Segmind Tiny SD (Portrait) Serverless APIs, Segmind offers fastest deployment for Tiny-Stable-Diffusion inferences

The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images

This model corresponds to the Stable Diffusion RPG checkpoint for detailed images at the cost of a super detailed prompt

This model corresponds to the Stable Diffusion Reliberate checkpoint for detailed images at the cost of a super detailed prompt

This model corresponds to the Stable Diffusion Realistic Vision checkpoint for detailed images at the cost of a super detailed prompt

The most versatile photorealistic model that blends various models to achieve the amazing realistic images

This model corresponds to the Stable Diffusion Paragon checkpoint for detailed images at the cost of a super detailed prompt

Stable Diffusion Outpainting can extend any image in any direction

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

This model corresponds to the Stable Diffusion Epic Realism checkpoint for detailed images at the cost of a super detailed prompt

This model corresponds to the Stable Diffusion Edge of Realism checkpoint for detailed images at the cost of a super detailed prompt

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

Dreamshaper excels in delivering high-quality, detailed images. It is fine-tuned to understand and interpret a diverse range of artistic styles and subjects.

This model corresponds to the Stable Diffusion Disney checkpoint for detailed images at the cost of a super detailed prompt

The most versatile photorealistic model that blends various models to achieve the amazing realistic space themed images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

This model corresponds to the Stable Diffusion Colorful checkpoint for detailed images at the cost of a super detailed prompt

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

The most versatile photorealistic model that blends various models to achieve the amazing realistic images.

Create beautiful and creative QR codes for your marketing campaigns.

Convert text to images with the distilled stable diffusion model by Segmind, Small-SD. Segmind Small SD Serverless APIs, Segmind offers fastest deployment for Small-Stable-Diffusion inferences.

Kandinsky inherits best practices from Dall-E 2 and Latent diffusion, while introducing some new ideas.

This model corresponds to the ControlNet conditioned on Soft Edge.

This model corresponds to the ControlNet conditioned on Scribble images.

This model corresponds to the ControlNet conditioned on Depth estimation.

This model corresponds to the ControlNet conditioned on Canny edges.

CodeFormer is a robust face restoration algorithm for old photos or AI-generated faces.

The Segment Anything Model (SAM) produces high quality object masks from input prompts such as points or boxes, and it can be used to generate masks for all objects in an image.

Take a picture/gif and replace the face in it with a face of your choice. You only need one image of the desired face. No dataset, no training

This model corresponds to the Stable Diffusion Revanimated checkpoint for detailed images at the cost of a super detailed prompt

This model removes the background image from any image

AI-Powered Image Super-Resolution, upscaling and Image enhancement producing stunning, high-quality results using artificial intelligence

This model corresponds to the ControlNet conditioned on Human Pose Estimation.